Create highly detailed anime style with Stable Diffusion.

Stable Diffusion fine-tuned on Midjourney v4 images.

kandinsky-2.2 inherits best practices from Dall-E 2 and Latent diffusion, while introducing some new ideas.

Practical face restoration algorithm for old photos or AI-generated faces.

Bringing Old Photos Back to Life

Colorization using a Generative Color Prior for Natural Images.

Image Restoration using SwinIR Transformer.

Get an approximate text prompt, with style, matching an image.

Text-to-audio generation with latent diffusion models.

Fill in masked parts of images with Stable Diffusion.

A large language model that's been fine-tuned on ChatGPT interactions

Edit images with human instructions

Answers questions about images

Generate detailed images from scribbled drawings

This model is designed to produce high-quality, highly detailed anime-style images with just a few prompts.

Upload a face image and get back a professional headshot

Upload a single face photo and see how that person would look like if they were: fat or fit

The fastest way to swap a face in a photo without using Photoshop

AI model of the new generation for improving image quality

Playground v2.5 is the state-of-the-art open-source model in aesthetic quality

Detailed, higher-resolution images created with Stable Diffusion.

A latent text-to-image diffusion model capable of generating photo-realistic images given any text input.

The CLIP Interrogator is a prompt engineering tool that combines OpenAI's CLIP and Salesforce's BLIP to optimise text prompts to match a given image. Use the resulting prompts with text-to-image models like Stable Diffusion to create cool art!

Swin2SR Transformer for Image Super-Resolution and Restoration.

Image Style Transfer with Text Condition.

Arbitrary Neural Style Transfer

Multi-language Text-To-Speech Audio Model

Text-Driven Manipulation of StyleGAN Imagery

One of the fastest text-to-image AI models from a well-known tech giant.

Add colours to old video footage. The processing speed varies greatly depending on the input data.

Generate any type of logo using text.

Generate photorealistic images using text.

Edit input image with the help of text prompts.

Generate a new image given any input text with DreamShaper V8

Remove images background

Generate music from a prompt or melody

Removes defocus blur in an image

Generate super realistic images using text prompts

Animate your personalized text-to-image diffusion models

Create your own variants of “This is Fine”

Make stickers with AI. Generates graphics with transparent backgrounds.

Turn a face into 3D, emoji, pixel art, video game, claymation or toy

Animate Stable Diffusion by interpolating between two prompts.

Extract the foreground of a video. The processing speed varies greatly depending on the input data and typically completes within 4 minutes.

Animating prompts with Stable Diffusion.

Text to Audio using iNstruction-Guided diffusiOn

SDXL is a text-to-image generative AI model developed by Stability AI that creates beautiful 1024x1024 images. It is the successor to Stable Diffusion.

A 70 billion parameter language model from Meta

Multilingual Text-To-Speech voice cloning model

Simplify QR code creation using ControlNet user-friendly neural interface

Transcribing audio files into text

Take a picture of your face and instantly get any profile picture you want. Only 1 photo, no training needed.

Create photos, paintings and avatars for anyone in any style within seconds.

Create logos with AI in just one minute

Automatically generate product advertising image

Generate videos by interpolating the latent space of Stable Diffusion using the Openjourney Model

Face swap on video in one click

Create song covers with any RVC v2 trained AI voice from audio files.

Audio-based Lip Synchronization for Talking Head Video

Controllable and Consistent Human Image Animation with 3D Parametric Guidance

Create videos from illustrated input images.